AIGC笔记--基于Stable Diffusion实现图片的inpainting

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猴君
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1--完整代码

SD_Inpainting

2--简单代码

import PIL import torch import numpy as np from PIL import Image from tqdm import tqdm import torchvision from diffusers import AutoencoderKL, UNet2DConditionModel, DDIMScheduler from transformers import CLIPTextModel, CLIPTokenizer  # 预处理mask def preprocess_mask(mask):     mask = mask.convert("L") # 转换为灰度图: L = R * 299/1000 + G * 587/1000+ B * 114/1000。     w, h = mask.size # 512, 512     w, h = map(lambda x: x - x % 32, (w, h))  # resize to integer multiple of 32     mask = mask.resize((w // 8, h // 8), resample = PIL.Image.NEAREST) # 64, 64     mask = np.array(mask).astype(np.float32) / 255.0 # 归一化 64, 64     mask = np.tile(mask, (4, 1, 1)) # 4, 64, 64     mask = mask[None].transpose(0, 1, 2, 3)     mask = 1 - mask  # repaint white, keep black # mask图中,mask的部分变为0     mask = torch.from_numpy(mask)     return mask  # 预处理image def preprocess(image):     w, h = image.size     w, h = map(lambda x: x - x % 32, (w, h))  # resize to integer multiple of 32     image = image.resize((w, h), resample=PIL.Image.LANCZOS)     image = np.array(image).astype(np.float32) / 255.0     image = image[None].transpose(0, 3, 1, 2)     image = torch.from_numpy(image)     return 2.0 * image - 1.0  if __name__ == "__main__":     model_id = "runwayml/stable-diffusion-v1-5" # online download     # model_id = "/mnt/dolphinfs/hdd_pool/docker/user/hadoop-waimai-aigc/liujinfu/All_test/test0714/huggingface.co/runwayml/stable-diffusion-v1-5" # local path      # 读取输入图像和输入mask     input_image = Image.open("./images/overture-creations-5sI6fQgYIuo.png").resize((512, 512))     input_mask = Image.open("./images/overture-creations-5sI6fQgYIuo_mask.png").resize((512, 512))      # 1. 加载autoencoder     vae = AutoencoderKL.from_pretrained(model_id, subfolder = "vae")      # 2. 加载tokenizer和text encoder      tokenizer = CLIPTokenizer.from_pretrained(model_id, subfolder = "tokenizer")     text_encoder = CLIPTextModel.from_pretrained(model_id, subfolder = "text_encoder")      # 3. 加载扩散模型UNet     unet = UNet2DConditionModel.from_pretrained(model_id, subfolder = "unet")      # 4. 定义noise scheduler     noise_scheduler = DDIMScheduler(         num_train_timesteps = 1000,         beta_start = 0.00085,         beta_end = 0.012,         beta_schedule = "scaled_linear",         clip_sample = False, # don't clip sample, the x0 in stable diffusion not in range [-1, 1]         set_alpha_to_one = False,     )      # 将模型复制到GPU上     device = "cuda"     vae.to(device, dtype = torch.float16)     text_encoder.to(device, dtype = torch.float16)     unet = unet.to(device, dtype = torch.float16)      # 设置prompt和超参数     prompt = "a mecha robot sitting on a bench"     negative_prompt = ""     strength = 0.75     guidance_scale = 7.5     batch_size = 1     num_inference_steps = 50     generator = torch.Generator(device).manual_seed(0)      with torch.no_grad():         # get prompt text_embeddings         text_input = tokenizer(prompt, padding = "max_length",              max_length = tokenizer.model_max_length,              truncation = True,              return_tensors = "pt")         text_embeddings = text_encoder(text_input.input_ids.to(device))[0]          # get unconditional text embeddings         max_length = text_input.input_ids.shape[-1]         uncond_input = tokenizer(             [negative_prompt] * batch_size, padding = "max_length", max_length = max_length, return_tensors = "pt"         )         uncond_embeddings = text_encoder(uncond_input.input_ids.to(device))[0]         # concat batch         text_embeddings = torch.cat([uncond_embeddings, text_embeddings])          # 设置采样步数         noise_scheduler.set_timesteps(num_inference_steps, device = device)          # 根据strength计算timesteps         init_timestep = min(int(num_inference_steps * strength), num_inference_steps)         t_start = max(num_inference_steps - init_timestep, 0)         timesteps = noise_scheduler.timesteps[t_start:]          # 预处理init_image         init_input = preprocess(input_image)         init_latents = vae.encode(init_input.to(device, dtype=torch.float16)).latent_dist.sample(generator)         init_latents = 0.18215 * init_latents         init_latents = torch.cat([init_latents] * batch_size, dim=0)         init_latents_orig = init_latents          # 处理mask         mask_image = preprocess_mask(input_mask)         mask_image = mask_image.to(device=device, dtype=init_latents.dtype)         mask = torch.cat([mask_image] * batch_size)                  # 给init_latents加噪音         noise = torch.randn(init_latents.shape, generator = generator, device = device, dtype = init_latents.dtype)         init_latents = noise_scheduler.add_noise(init_latents, noise, timesteps[:1])         latents = init_latents # 作为初始latents          # Do denoise steps         for t in tqdm(timesteps):             # 这里latens扩展2份,是为了同时计算unconditional prediction             latent_model_input = torch.cat([latents] * 2)             latent_model_input = noise_scheduler.scale_model_input(latent_model_input, t) # for DDIM, do nothing              # 预测噪音             noise_pred = unet(latent_model_input, t, encoder_hidden_states=text_embeddings).sample              # Classifier Free Guidance             noise_pred_uncond, noise_pred_text = noise_pred.chunk(2)             noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond)              # x_t -> x_t-1             latents = noise_scheduler.step(noise_pred, t, latents).prev_sample                          # 将unmask区域替换原始图像的nosiy latents             init_latents_proper = noise_scheduler.add_noise(init_latents_orig, noise, torch.tensor([t]))             # mask的部分数值为0             # 因此init_latents_proper * mask为保留原始latents(不mask)             # 而latents * (1 - mask)为用生成的latents替换mask的部分             latents = (init_latents_proper * mask) + (latents * (1 - mask))           # 注意要对latents进行scale         latents = 1 / 0.18215 * latents         image = vae.decode(latents).sample                  # 转成pillow         img = (image / 2 + 0.5).clamp(0, 1).detach().cpu()         img = torchvision.transforms.ToPILImage()(img.squeeze())         img.save("./outputs/output.png")         print("All Done!")

运行结果:

3--基于Diffuser进行调用

import torch import torchvision from PIL import Image from diffusers import StableDiffusionInpaintPipelineLegacy  if __name__ == "__main__":     # load inpainting pipeline     model_id = "runwayml/stable-diffusion-v1-5"     # model_id = "/mnt/dolphinfs/hdd_pool/docker/user/hadoop-waimai-aigc/liujinfu/All_test/test0714/huggingface.co/runwayml/stable-diffusion-v1-5" # local path     pipe = StableDiffusionInpaintPipelineLegacy.from_pretrained(model_id, torch_dtype = torch.float16).to("cuda")      # load input image and input mask     input_image = Image.open("./images/overture-creations-5sI6fQgYIuo.png").resize((512, 512))     input_mask = Image.open("./images/overture-creations-5sI6fQgYIuo_mask.png").resize((512, 512))      # run inference     prompt = ["a mecha robot sitting on a bench", "a cat sitting on a bench"]     generator = torch.Generator("cuda").manual_seed(0)     with torch.autocast("cuda"):         images = pipe(             prompt = prompt,             image = input_image,             mask_image = input_mask,             num_inference_steps = 50,             strength = 0.75,             guidance_scale = 7.5,             num_images_per_prompt = 1,             generator = generator         ).images      # 转成pillow     for idx, image in enumerate(images):         image.save("./outputs/output_{:d}.png".format(idx))     print("All Done!")

运行结果:

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